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. 2018 Mar 22;372(3):1597–1629. doi: 10.1007/s00208-018-1672-1

Hyperbolic systems with non-diagonalisable principal part and variable multiplicities, I: well-posedness

Claudia Garetto 1, Christian Jäh 1, Michael Ruzhansky 2,
PMCID: PMC6411233  PMID: 30930490

Abstract

In this paper we analyse the well-posedness of the Cauchy problem for a rather general class of hyperbolic systems with space-time dependent coefficients and with multiple characteristics of variable multiplicity. First, we establish a well-posedness result in anisotropic Sobolev spaces for systems with upper triangular principal part under interesting natural conditions on the orders of lower order terms below the diagonal. Namely, the terms below the diagonal at a distance k to it must be of order -k. This setting also allows for the Jordan block structure in the system. Second, we give conditions for the Schur type triangularisation of general systems with variable coefficients for reducing them to the form with an upper triangular principal part for which the first result can be applied. We give explicit details for the appearing conditions and constructions for 2×2 and 3×3 systems, complemented by several examples.

Mathematics Subject Classification: 35L45 (primary), 46E35 (secondary)

Introduction

The main aim of this paper is to consider the Cauchy problem for hyperbolic systems

Dtu=A(t,x,Dx)u+B(t,x,Dx)u+f(t,x),(t,x)[0,T]×Rn,ut=0=u0,xRn, 1

with the usual notation Dt=-it and Dx=-ix. Here, we assume that A(t,x,Dx)=[aij(t,x,Dx)]i,j=1m is an m×m matrix of pseudo-differential operators of order 1, i.e. aijC([0,T],Ψ1,01(Rn)) with possibly complex valued symbols. In the first part of the paper we will also assume that

A(t,x,Dx)=Λ(t,x,Dx)+N(t,x,Dx), 2

with real-valued symbols in

Λ(t,x,Dx)=diag(λ1(t,x,Dx),λ2(t,x,Dx),,λm(t,x,Dx)),

and

N(t,x,Dx)=0a12(t,x,Dx)a13(t,x,Dx)a1m(t,x,Dx)00a23(t,x,Dx)a2m(t,x,Dx)000am-1m(t,x,Dx)0000.

Finally, we assume that

B(t,x,Dx)=[bij(t,x,Dx)]i,j=1m,bijC([0,T],Ψ1,00(Rn)),

is an m×m matrix of pseudo-differential operators of order 0 with possibly complex valued symbols. We can take any n1 and we can assume that m2 since in the case m=1 there are no multiplicities and thus much more is known. It is also well-known that even if all the coefficients in A and B depend only on time, due to multiplicities, the best one can hope for is the well-posedness of the Cauchy problem (1) in suitable classes of Gevrey spaces. Thus, the main questions that we address in this paper are:

  1. Under what structural conditions on the zero order part B(t,x,Dx) is the Cauchy problem (1) well-posed in C or, even better, in suitable scales of Sobolev spaces?

  2. Under what conditions on the general matrix A(t,x,Dx) of first order pseudo-differential operators can we reduce it (microlocally) to another system with A satisfying the upper triangular condition (2)?

Note that this paper is part of a wider analysis of hyperbolic systems with multiplicities. Here we investigate the well-posedness of these systems. In the second part of this paper we plan to carry out the microlocal analysis of their solutions.

In the case of 2×2 systems the questions above have been analysed with the answer to (Q1) given by the following theorem:

Theorem A

([27, Theorem 7.2]) Let m=2. Suppose that the pseudo-differential operator b21 is of order not greater than -1. Then the Cauchy problem (1) is well-posed in C. Moreover, it is well-posed in the anisotropic Sobolev space Hs1(Rn)Hs2(Rn) provided s2-s11. In that case the solution satisfies the following estimates:

u1(t,·)Hs+u2(t,·)Hs+1cectu10Hs+u20Hs+1,0tT,

for uj0Hcomps+j-1(Rn), j=1,2, with c>0 depending on s, T, and the support of the initial data.

The case of systems of general size but for coefficients depending only on t and for n=1 was also considered. More precisely, in [26] the authors considered the Cauchy problem

Dtu=A(t)Dxu+B(t)Dxu+f(t,x),(t,x)[0,T]×R,ut=0=0,xR, 3

with A(t)=[aij(t)]i,j=1mC([0,T])m×m in the form

A(t)=Λ(t)+N(t),

similar to (2). They showed the following result in the absence of lower order terms and for zero Cauchy data:

Theorem B

([26, Proposition 1]) Let B(t)0 and let sR. Then the Cauchy problem (3) is C-well-posed. Moreover, there exist r1,..., rm-1[0,1] such that for every fC(R,(Hs(R))m) identically 0 at t=0 it admits a unique solution uC(R,(S(R))m) satisfying

umC(R,Hs(R)),um-jC(R,Hs-r1--rj-1(R)),

for j=1,,m-1, and identically 0 at t=0. In particular, if λj(t)λk(t), tR, 1j<km, no loss of anisotropic regularity appears.

The case of (microlocally) diagonalisable systems of any order with fully variable coefficients was considered by Rozenblum [41] under the condition of transversality of the intersecting characteristics. Also allowing the variable multiplicities, this transversality condition was later removed in [32, 33] with sharp Lp-estimates for solutions, with further applications to the spectral asymptotics of the corresponding elliptic systems.

Before stating our main results and collecting some necessary basic notions we give a brief overview of the state of the art for hyperbolic equations and systems. We have a complete understanding of strictly hyperbolic systems, i.e., systems without multiplicities, with C-coefficients. This starts with the groundbreaking work of Lax [35] and Hörmander [28] and heavily relies on the modern theory of Fourier integral operators (FIO). Well-posedness is here obtained in the space of distributions D. There are also well-posedness results for less regular coefficients with respect to t. For instance, well-posedness with loss of derivatives has been obtained by Colombini and Lerner [9] for second order strictly hyperbolic equations with Log-Lipschitz coefficients with respect to t and smooth in x. It is possible to further drop the regularity in t (for instance Hölder), however, this has to be balanced by stronger regularity in x (Gevrey) and leads to more specific (Gevrey) well-posedness results (see [3, 31] and references therein). Paradifferential techniques have been recently used for this kind of strictly hyperbolic equations by Colombini et al. [6, 7].

The analysis of hyperbolic equations with multiplicities (weakly hyperbolic) has started with the seminal paper by Colombini et al. [5] in the case of coefficients depending only on time. Profound difficulties in such analysis have been exhibited by Colombini et al. [4, 8] showing that even the second order wave equation in R with smooth time-dependent propagation speed (but with multiplicity) and smooth Cauchy data need not be well-posed in D. However, they turn out to be well-posed in suitable Gevrey classes or spaces of ultradistributions. In the last decades many results were obtained for weakly hyperbolic equations with t-dependent coefficients ([3, 11, 16, 1820, 34], to quote only very few). More recently, advances in the theory of weakly hyperbolic systems with t-dependent coefficients have been obtained for systems of any size in presence of multiplicities with regular or low regular (Hölder) coefficients [16, 22, 23]. In addition, in [17] precise conditions on the lower order terms (Levi conditions) have been formulated to guarantee Gevrey and ultradistributional well-posedness. Previously very few results were known in the field for systems of a certain size (2×2, 3×3) [12, 13] or of a certain form (for instance without lower order terms or with principal part of a certain form) [44].

Weakly hyperbolic equations with x-dependent coefficients were considered for the first time in the celebrated paper by Bronshtein [2]. As shown already in some earlier works by Ivrii, the corresponding Cauchy problem is well-posed under “almost analytic regularity”, namely, if the coefficients and initial data are in suitable Gevrey classes. Bronshtein’s result was extended to (tx)-dependent scalar equations by Ohya and Tarama [38] and to systems by Kajitani and Yuzawa [31]. The regularity assumptions are always quite strong with respect to x (Gevrey) and not below Hölder in t. See also [10, 37]. Geometrical and microlocal analytic approaches are known for equations or systems under specific assumptions on the characteristics and/or lower order terms. See [29, 30, 33, 36, 39], to quote only a few. Time-dependent coefficients of low regularity (distributional) have been considered in [21].

In this paper we will be interested in the case of coefficients depending on both t and x and we will make use of the usual definitions of symbol classes. We say that a (possibly) complex valued function a=a(x,ξ)C(Rn×Rn) belongs to S1,0m(Rn×Rn) if there exist constants Cα,β>0 such that

α,βN0n:|xαξβa(x,ξ)|Cα,βξm-|β|(x,ξ)Rn×Rn.

The set of pseudo-differential operators associated to the symbols in S1,0m(Rn×Rn) is denoted by Ψ1,0m(Rn×Rn).

If there is no question about the domain under consideration, we will abbreviate the symbol- and operator-classes by S1,0m and Ψ1,0m, respectively, or simply by Sm and Ψm.

We also denote by C([0,T],S1,0m(Rn×Rn)) the space of all symbols a(t,x,ξ)S1,0m(Rn×Rn) which are continuous with respect to t. The set of operators associated to the symbols in C([0,T],S1,0m(Rn×Rn)) is denoted by C([0,T],Ψ1,0m(Rn×Rn)).

Again, if there is no question about the domain under consideration, we will abbreviate the symbol- and operator-classes by CS1,0m and CΨ1,0m, respectively, or simply by CSm and CΨm.

Let us give our main result concerning the first question (Q1) for the systems with the principal part A satisfying the upper triangular condition (2). Here, fk, uk and uk0, for k=1,,m, stand for the components of the vectors f, u and u0, respectively.

Theorem 1

Let n1, m2, and let

Dtu=A(t,x,Dx)u+B(t,x,Dx)u+f(t,x),(t,x)[0,T]×Rn,ut=0=u0(x),xRn, 4

where A(t,x,Dx)(CS1)m×m is an upper-triangular matrix of pseudo-differential operators of order 1 in the form (2), and B(t,x,Dx)(CS0)m×m is a matrix of pseudo-differential operators of order 0, continuous with respect to t. Hence, if

the lower order termsbijbelong toC([0,T],Ψj-i)fori>j,

uk0Hs+k-1(Rn) and fkC([0,T],Hs+k-1) for k=1,,m, then (4) has a unique anisotropic Sobolev solution u, i.e., ukC([0,T],Hs+k-1) for k=1,,m.

Remark 1

As stated earlier, we allow A and B to have complex valued symbols as long as the symbols of Λ in (2), i.e. the eigenvalues of A(t,x,ξ), are real valued.

The main condition of Theorem 1 for the Sobolev well-posedness is that the pseudo-differential operator bij below the diagonal (i.e. for i>j) must be of order j-i. In other words, the terms below the diagonal at a distance k to it must be of order -k.

In solving the Cauchy problem (4) we will make use of Fourier integral operators depending on the parameter t[0,T]. Namely, we will work with operators of the type

0tRneiφ(t,s,x,ξ)a(t,s,x,ξ)g^(s,ξ)dξds

where φ is the solution of a certain eikonal equation and the symbol a is determined via asymptotic expansion and transport equations. In Sect. 2.1 we will recall some well-known Sobolev estimates for this type of operators.

In Sect. 2 we will prove Theorem 1 after we explain its idea in the cases of m=2 and m=3.

Consequently, in Sect. 3 we give an answer to the second question (Q2) above in the form of a suitable variable coefficients extension of the Schur triangularisation. For constant matrices such a procedure is well known (see e.g. [1, Theorem 5.4.1]).

Theorem C

(Schur’s triagularisation theorem) Given a (constant) m×m matrix A with eigenvalues λ1,,λm in any prescribed order, there is a unitary m×m matrix T such that R=T-1AT is upper triangular with the diagonal elements rii=λi. Furthermore, if the entries of A and its eigenvalues are all real, T may be chosen to be real orthogonal.

It follows that R can be written as D+N, where D=diag(λ1,,λm) and N is a nilpotent upper triangular matrix.

If the matrix A depends on one or several parameters, namely A=A(t,x,ξ), the situation becomes less clear and it is difficult to give a complete description, in particular together with a prescribed regularity of the involved transformation matrices. The regularity of the matrix A and the desire to maintain it through the transformation puts already constrains on the matrix as, in general, the eigenvalues can only be expected to be Lipschitz continuous in the parameters even if all the entries depend smoothly on the parameters (see, e.g., [2, 40] and the references therein). In the sequel, we will present some sufficient conditions to ensure the existence of an upper triangularisation for A(t,x,ξ) which respects its regularity. For example, it will apply to the case when A is a matrix of first order symbols continuous with respect to t, i.e., A(t,x,ξ)(CS1)m×m.

Our main result for this part of the problem is the following theorem.

Theorem 2

Let A(t,x,ξ)(CS1)m×m, be a m×m-matrix with eigenvalues λ1,,λmCS1, and let h1,,hm-1(CS0)m be the corresponding eigenvectors. Suppose that for e1=[1,0,,0]TRm-i+1 the condition

h(i)(t,x,ξ)|e10,(t,x,ξ)[0,T]×Rn×Rn 5

holds for all i=1,,m-1, with the notation for h(i) explained in (37). Then, there exists a matrix-valued symbol T(t,x,ξ)(CS0)m×m, invertible for (t,x,ξ)[0,T]×Rn×{|ξ|M} with T-1(t,x,ξ)(CS0)m×m, such that

T-1(t,x,ξ)A(t,x,ξ)T(t,x,ξ)=Λ(t,x,ξ)+N(t,x,ξ)

for all (t,x,ξ)[0,T]×Rn×{|ξ|M}, where

Λ(t,x,ξ)=diag(λ1(t,x,ξ),λ2(t,x,ξ),,λm(t,x,ξ))

and

N(t,x,ξ)=0N12(t,x,ξ)N13(t,x,ξ)N1m(t,x,ξ)00N23(t,x,ξ)N2m(t,x,ξ)000Nm-1m(t,x,ξ)0000,

and N is a nilpotent matrix with entries in CS1.

Furthermore, there is an expression for the matrix symbol T which will be given in Theorem 6. Also, the assumption (5) can be relaxed, see Remark 6. In Sect. 3 we will prove this result as well as describe the procedure how to obtain the desired upper triangular form. Moreover, we work out in detail the cases of m=2 and m=3 clarifying this Schur triangualisation procedure and give a number of examples.

The results and techniques of this paper are a natural outgrowth of the paper [27] where the case m=2 was considered and to which the results of the present paper reduce in the case of 2×2 systems. It is with great sorrow that we remember the untimely departure of our colleague and friend Todor Gramchev who was the inspiration for both [27] and the present paper.

Well-posedness in anisotropic Sobolev spaces

This section is devoted to proving the well-posedness of the Cauchy problem (1). For the reader’s convenience we first give a detailed proof in the cases m=2 and m=3. This will inspire us in proving Theorem 1. We note that the case m=2 has been studied in [27] and we will briefly review its derivartion. However, first we collect a few results about Fourier integral operators that we will need in the sequel.

Auxiliary remarks

In solving the Cauchy problem (1), we will deal with solutions of certain scalar pseudo-differential equations. For each characteristic λj of A, we will be denoting by Gj0θ the solution to

Dtw=λj(t,x,Dx)w+bjj(t,x,Dx)w,w(0,x)=θ(x),

and by Gjg the solution to

Dtw=λj(t,x,Dx)w+bjj(t,x,Dx)w+g(t,x),w(0,x)=0.

The operators Gj0 and Gj can be microlocally represented by Fourier integral operators

Gj0θ(t,x)=Rneiφj(t,x,ξ)aj(t,x,ξ)θ^(ξ)dξ 6

and

Gjg(t,x)=0tRneiφj(t,s,x,ξ)Aj(t,s,x,ξ)g^(s,ξ)dξds,

with φj(t,s,x,ξ) solving the eikonal equation

tφj=λj(t,x,xφj),φj(s,s,x,ξ)=x·ξ,

and with the notation

φj(t,x,ξ)=φj(t,0,x,ξ).

Here we also have the amplitudes Aj,-k(t,s,x,ξ) of order -k, k N, giving Ajk=0Aj,-k, and they satisfy the transport equations with initial data at t=s, and we have aj(t,x,ξ)=Aj(t,0,x,ξ).

If ajSm, i.e. if the amplitude aj in (6) is a symbol of order m, we will write Gj0I1,0m. However, in the above construction of propagators for hyperbolic equations, we have ajS0, so that Gj0I1,00.

By I1,0m, we denote the class of Fourier integral operators with amplitudes in S1,0m. For further information, the reader may consult [15, 42, 43] and the references therein.

With that, we can record the following estimate:

Lemma 1

For any σR, for sufficiently small t, we have

Gj0θ(t)HσCA,σ,u0θHσ,Gjg(t)HσCA,σtgLsHxσ.

This statement follows from the continuity of λj,φj,aj,Aj with respect to t and from the Hσ-boundedness of non-degenerate Fourier integral operators, see e.g. [15] (there are also surveys on such questions [42, 43]). It is important to note that the constant for the estimate for Gj does not depend on the initial data of the Cauchy problem; see also Remark 2.

The case m=2

To motivate the higher order cases, here we review the construction for 2×2 systems adapting it for the subsequent higher order arguments. Hence, in this subsection we follow the proof in [27]. Thus, we consider the system

Dtu=A(t,x,Dx)u+B(t,x,Dx)u+f(t,x),(t,x)[0,T]×Rn,ut=0=u0,xRn, 7

where u0(x)=[u10(x),u20(x)]T, f(t,x)=[f1(t,x),f2(t,x)]T, and with the operators A(t,x,Dx) and B(t,x,Dx) given by

A(t,x,Dx)=λ1(t,x,Dx)a12(t,x,Dx)0λ2(t,x,Dx) 8

and

B(t,x,Dx)=b11(t,x,Dx)b12(t,x,Dx)b21(t,x,Dx)b22(t,x,Dx).

We suppose that all entries of A(t,x,Dx) belong to CΨ1,01 and all entries of B(t,x,Dx) belong to CΨ1,00. By using the operators Gj0 and Gj introduced in Sect. 2.1, we can reformulate the Eq. (7) as

u1=U10+G1((a12+b12)u2), 9
u2=U20+G2(b21u1), 10

where

Uj0=Gj0uj0+Gj(fj),j=1,2. 11

Plugging (10) in (9), we obtain

u1=U~10+G1(a12G2(b21u1))+G1(b12G2(b21u1)), 12

where

U~10=G10u10+G1(f1)+G1((a12+b12)U20). 13

Using the rules of composition of Fourier integral operators, see e.g. [15], and by Lemma 1, we get that the operator G1a12G2b21 in (12) acts continuously on Hs if it is of order 0. Since a12CΨ1,01 we therefore need to assume that b21CΨ1,0-1.

The operator G1b12G2b21 belongs to CI1,0-1 since b21CΨ1,0-1 and b12CΨ1,00.

We now introduce the following scale of Banach spaces Xs(t):=C([0,t],Hs), t[0,T], equipped with the norm

uXs(t)=supτ[0,t]u(τ,·)Hs.

Let

G10u1:=G1(a12G2(b21u1))+G1(b12G2(b21u1)).

It follows that (12) can be written as

u1=U~10+G10u1.

By composition of Fourier integral operators and Lemma 1 we have that the 0-order Fourier integral operator G10 maps C([0,T],Hs) continuously into itself and for small time interval it is a contraction, in the sense that there exists T[0,T] such that

G10(u-v)Xs(T)CA,sTu-vXs(T),

with CA,sT<1. Banach’s fixed point theorem ensures the existence of a unique fixed point u1 for the map G10. Hence, by assuming that the initial data U~10 belongs to C([0,T],Hs) we conclude that there exists a unique u1C([0,T],Hs) solving (12). Note that the same argument proves that the operator I-G10 is invertible on a sufficiently small interval in t since G10=I at t=0. From formula (13) it is clear that in order to get U~10 to belong to C([0,T],Hs) we need to assume that U20Hs+1. Finally, we get u2 by substitution of u1 in (10).

Remark 2

Note that the constant T depends only on A and s. Thus, the argument above can be iterated by taking u(T,x) as new initial data. In this way one can cover an arbitrary finite interval [0, T] and obtain a solution in C([0,T],Hs)×C([0,T],Hs+1).

Remark 3

Since a12(t,x,Dx) is a first order operator combining (11) with (13) we easily see that in order to get Sobolev well-posedness of order s we need to take initial data u10 and u20 in Hs and Hs+1, respectively, and right hand-side functions f1 and f2 in C([0,T],Hs) and C([0,T],Hs+1), respectively.

We have therefore proved the following theorem stated for the first time in [27, Theorem 7.2].

Theorem 3

Consider the Cauchy problem (7), with the 2×2 matrices

A(t,x,Dx)(CS1)2×2andB(t,x,Dx)(CS0)2×2,

where A is of the form (2). Assume that b21C([0,T],Ψ1,0-1), the right hand-side functions f1 and f2 belong to C([0,T],Hs) and C([0,T],Hs+1), respectively, and the initial data u10 and u20 belong to Hs and Hs+1, respectively. Then, (7) has a unique solution in C([0,T],Hs)×C([0,T],Hs+1). More generally it is well-posed in the anisotropic Sobolev space C([0,T],Hs1)×C([0,T],Hs2), provided s2-s1=1.

Remark 4

It was also shown in [27] that the solution u satisfies the estimate

u1(t,·)Hs+u2(t,·)Hs+1cectu10Hs+u20Hs+1,0tT,

for uj0Hcomps+j-1, j=1,2 with c>0 depending on s, T, and the support of the initial data. Since well-posedness is obtained for any Sobolev order s it follows that the Cauchy problem (7) is also C well-posed.

The case m=3

In this section we will extend the construction to the case of 3×3 systems. In the argument there is an additional substitution and a fixed point argument step compared to the case m=2. The advantage of giving the case of m=3 here is that we can make the argument more concrete compared to the more abstract construction in the general case that will be given in the following section. Thus, let

Dtu=A(t,x,Dx)u+B(t,x,Dx)u+f(t,x),(t,x)[0,T]×Rn,ut=0=u0,xRn, 14

where u0(x)=[u10(x),u20(x),u30(x)]T, f(t,x)=[f1(t,x),f2(t,x),f3(t,x)]T, A(t,x,Dx) is defined by the matrix

λ1(t,x,Dx)a12(t,x,Dx)a13(t,x,Dx)0λ2(t,x,Dx)a23(t,x,Dx)00λ3(t,x,Dx), 15

and

B(t,x,Dx)=b11(t,x,Dx)b12(t,x,Dx)b13(t,x,Dx)b21(t,x,Dx)b22(t,x,Dx)b23(t,x,Dx)b31(t,x,Dx)b32(t,x,Dx)b33(t,x,Dx).

We assume that all the entries of A(t,x,Dx) and B(t,x,ξ) belong to CΨ1,01 and CΨ1,00, respectively. Using the notations introduced earlier, we can write

u3(t,x)=U30+G3(b31u1)+G3(b32u2),u2(t,x)=U20+G2((a23+b23)u3)+G2(b21u1),u1(t,x)=U10+G1((a12+b12)u2)+G1((a13+b13)u3), 16

where

Uj0(t,x)=Gj0(uj0)+Gj(fj),j=1,2,3. 17

Now, we plug u3 into u1 and u2 in formula (16) and, thus, obtain

u2(t,x)=U~20+G2(b21u1)+G2((a23+b23)G3(b31u1))+G2((a23+b23)G3(b32u2)),u1(t,x)=U~10+G1((a13+b13)G3(b31u1))+G1((a13+b13)G3(b32u2))+G1((a12+b12)u2), 18

where

U~j0=Uj0+Gj((aj3+bj3)(t,x,Dx)U30),j=1,2.

We introduce the operator G20 by setting

G20u2:=G2((a23+b23)G3(b32u2)) 19

and in analogy with the case m=2 we define

L2u2:=u2-G20u2.

By Lemma 1 we have that for any s, G20 has the operator norm in Hs strictly less than 1 on a sufficiently small interval [0,T], so L2 is a perturbation of the identity operator. By the Neumann series it follows that L2 is invertible as a continuous operator from C([0,T],Hs) to C([0,T],Hs). Noting now that

u2-G20u2=L2u2=U~20+G2(b21u1)+G2((a23+b23)G3(b31u1)),

we have that

u2(t,x)=L2-1U~20+L2-1G2((a23+b23)G3(b31u1))+L2-1G2(b21u1).

Since this expression depends only on u1, we can plug it into the formula for u1 in (18) and obtain

u1(t,x)=U~10+G1((a13+b13)G3(b31u1))++G1((a13+b13)G3(b32u2))+G1((a12+b12)u2)=U~10+G1((a13+b13)G3(b31u1))++G1((a13+b13)G3(b32(L2-1U~20))+G1((a13+b13)G3(b32L2-1G2((a23+b23)G3(b31u1))))+G1((a13+b13)G3(b32(L2-1G2(b21u1)))+G1((a12+b12)L2-1U~20)+G1((a12+b12)L2-1G2((a23+b23)G3(b31u1)))+G1((a12+b12)L2-1G2(b21u1)).

By collecting now the terms with order 0 we can simplify the previous formula as follows:

u1(t,x)=U~10+G1(a13G3(b31u1))+G1(a13G3(b32(L2-1U~20)))+G1(a13G3b32L2-1G2(a23G3(b31u1)))+G1(a13G3b32L2-1G2(b23G3(b31u1)))+G1(a13G3b32(L2-1G2(b21u1)))+G1(b13G3(b32L2-1G2(a23G3(b31u1))))+G1(a12L2-1U~20)+G1(a12L2-1G2(a23G3(b31u1)))+G1(a12L2-1G2(b23)G3(b31u1)))+G1(b12L2-1G2((a23G3(b31u1)))+G1(a12L2-1G2(b21u1))+l.o.t.

Looking at the terms

G1(a13G3(b32(L2-1U~20))),G1(a12L2-1G2(b21u1)),G1(a12L2-1G2(a23G3(b31u1)))

and keeping in mind that in order to get the right Sobolev regularity we need to have operators of order 0, we deduce that b21 and b32 must have order -1 while b31 must have order -2. Considering now the initial data

U~j0=Uj0+Gj((aj3+bj3)U30),j=1,2,

by using (17) we obtain

U~j0=Uj0+Gj((aj3+bj3)(G30(u30)+G3(f3))),j=1,2.

Combining these formulas with an analysis of the term G1(a12L2-1U~20) we deduce that U~20 must belong to Hs+1. This implies U20Hs+1 and U30Hs+2. Concluding, similarly to the case m=2, that is by the Banach fixed point theorem argument on u1 and substitution in u2 and u3, we get anisotropic Sobolev well-posedness by assuming u10 and f1 in Hs, u20 and f2 in Hs+1, and u30 and f3 in Hs+2. This well-posedness is obtained by means of one invertible operator L2, and in analogy with case m=2 the well-posedness can be extended to the whole interval [0, T] by an iterated argument. This proves Theorem 1 in the case m=3.

The general case

We are now ready to prove the main result of our paper in the general case of an upper-triangular m×m matrix, i.e, a matrix A of the type

λ1(t,x,Dx)a12(t,x,Dx)a1m(t,x,Dx)0λ2(t,x,Dx)a2m(t,x,Dx)00λm-1(t,x,Dx)am-1m(t,x,Dx)00λm(t,x,Dx).

For the convenience of the reader we recall here the statement of Theorem 1.

Theorem 1

Let

Dtu=A(t,x,Dx)u+B(t,x,Dx)u+f(t,x),(t,x)[0,T]×Rn,ut=0=u0(x),xRn, 20

where A(t,x,Dx) is an upper-triangular matrix of pseudo-differential operators of order 1 in the form (2), and B(t,x,Dx) is a matrix of pseudo-differential operators of order 0, continuous with respect to t. Hence, if the lower order terms bij belong to C([0,T],Ψj-i) for i>j, uk0Hs+k-1 and fkC([0,T],Hs+k-1) for k=1,,m then (20) has a unique anisotropic Sobolev solution u, i.e., ukC([0,T],Hs+k-1) for k=1,,m.

Proof

Making use of the notations introduced earlier we can write the components of the solution u as

ui(t,x)=Ui0+Gij>imaij(t,x,Dx)uj+Gij=1jimbij(t,x,Dx)uj=Ui0+j<iGi(bij(t,x,Dx)uj)+i<jmGi((aij+bij)(t,x,Dx)uj), 21

where

Ui0=Gi0uj0+Gi(fi),

and Gi,Gi0 are Fourier integral operator of order 0 for i=1,,m. Note that from the fact that bij is a symbol of order 0 for every ij and, in particular, of order j-i for j<i we obtain that the operator Gi(bij) is of order j-i for j<i, while Gi(aij+bij) is, in general, of order 1. To simplify the argument we introduce the notations Gi,jj-i and Gi,j1 for the operators Gi(bij) and Gi(aij+bij), respectively. Here the superscript stands to remind us of the order of the operator. Hence,

ui=Ui0+j<iGi,jj-i(uj)+i<jmGi,j1(uj),

for i=1,,m. By begin by substituting

um=Um0+j<mGm,jj-m(uj),

into

um-1=Um-10+j<m-1Gm-1,jj-m+1(uj)+Gm-1,m1(um).

We get

um-1=Um-10+j<m-1Gm-1,jj-m+1(uj)+Gm-1,m1Um0+j<mGm-1,m1Gm,jj-m(uj)=(Um-10+Gm-1,m1Um0)+j<m-1(Gm-1,jj-m+1(uj)+Gm-1,m1Gm,jj-m(uj))+Gm-1,m1Gm,m-1-1um-1.

Note that it is enough to assume Um0Hs+1 and Um-10Hs to obtain Um-10+Gm-1,m1Um0Hs. Since all the operators above are of order 0 we conclude that the operator

Lm-1=I-Gm-1,m1Gm,m-1-1:=I-Gm-10

is invertible on a sufficiently small interval [0, T] and, therefore,

um-1-Gm-1,m1Gm,m-1-1um-1=(Um-10+Gm-1,m1Um0)+j<m-1(Gm-1,jj-m+1(uj)+Gm-1,m1Gm,jj-m(uj)), 22

yields

um-1=Lm-1-1U~m-10+Lm-1-1j<m-1G~m-1j-m+1uj, 23

with U~m-10 and G~m-1j-m+1 defined by the right-hand side of (22). We now substitute um and um-1 into um-2 making use of (23). We obtain

um-2=Um-20+j<m-2Gm-2,jj-m+2(uj)+Gm-2,m-11(um-1)+Gm-2,m1(um)=Um-20+j<m-2Gm-2,jj-m+2(uj)+Gm-2,m-11Lm-1-1U~m-10+Gm-2,m-11Lm-1-1j<m-2G~m-1j-m+1uj+Gm-2,m-11Lm-1-1G~m-1-1um-2+Gm-2,m1Um0+Gm-2,m1j<m-2Gm,jj-m(uj)+Gm-2,m1Gm,m-2-2um-2+Gm-2,m1Gm,m-1-1Lm-1-1U~m-10+Gm-2,m1Gm,m-1-1Lm-1-1j<m-2G~m-1j-m+1uj+Gm-2,m1Gm,m-1-1Lm-1-1G~m-1-1um-2. 24

We set

U~m-20=Um-20+Gm-2,m-11Lm-1-1U~m-10+Gm-2,m1Um0+Gm-2,m1Gm,m-1-1Lm-1-1U~m-10. 25

The operators Gm-2,m-11Lm-1-1 and Gm-2,m1 in (25) are of order 1. Keeping in mind that we already assumed Um0Hs+1 and Um-10Hs, in order to obtain Sobolev order s the initial data Um0, Um-10 and Um-20 must belong to Hs+2, Hs+1 and Hs, respectively. Thus,

um-2=U~m-20+Gm-20um-2+j<m-2G~m-2j-m+2uj, 26

where Gm-20 is a zero order operator defined by

Gm-20um-2=Gm-2,m-11Lm-1-1G~m-1-1um-2+Gm-2,m1Gm,m-2-2um-2+Gm-2,m1Gm,m-1-1Lm-1-1G~m-1-1um-2,

and the last summand in (26) is obtained by collecting all the operators acting on uj with j<m-2 in (24). Since the norm of Gm-20 can be taken strictly less than one in a sufficiently small interval [0, T] we have that the operator

Lm-2=I-Gm-20

is invertible and, therefore,

um-2=Lm-2-1U~m-20+j<m-2Lm-2-1G~m-2j-m+2uj. 27

Note that U~m-20Hs if Um0Hs+2, Um-10Hs+1 and Um-20Hs. By iterating the same procedure we deduce that

uk=U~k0+Gk0uk+j<kG~kj-kuj, 28

where U~k0 depends on Uk0, Uj0 and U~j0 with j>k and Gk0 is a zero order operator defined by using invertible operators Lm-1, Lm-2,..., Lk. In addition, we obtain U~k0Hs since Um0Hs+m-k, Um-10Hs+m-k-1,,Uk0Hs. It follows that for k=2 we have

u2=U~20+G20u2+G~2-1u1,

where the operator G20 is of zero operator and defined by invertible operators Lm-1,Lm-2,,L2, G~2-1 is of order -1, and U~20Hs since Um0Hs+m-2, Um-10Hs+m-3,,U20Hs. Hence, by inverting the operator L2=I-G20 on a sufficiently small interval [0, T] we have

u2=L2-1U~20+L2-1G~2-1u1.

Now by substitution of u2,u3,,um in the equation of u1 we arrive at the formula (28) with k=1, i.e.,

u1=U~10+G10u1,

where U~10Hs since Um0Hs+m-1, Um-10Hs+m-2,,U20Hs+1,U10Hs. Concluding, by the Banach fix point argument we prove that there exists a unique u1C([0,T],Hs) solving the equation above with the given initial conditions. By substitution in the equations for u2,,um-1,um we arrive at the desired Sobolev well-posedness with ukC([0,T],Hs+k-1) for k=2,,m. Note that, since the sufficiently small interval [0, T] where we get well-posedness does not depend on the initial data, by a standard iteration argument we can achieve well-posedness on any bounded interval [0, T] as stated in the theorem.

Schur decomposition of m×m matrices

In this section we investigate how to reduce an m×m matrix to the upper triangular form. We recall that such decomposition is well-known for constant matrices and goes under the name of Schur’s triangularisation, with its statement given in Theorem C.

One of the difficulties when dealing with variable multiplicities is the loss of regularity in the parameters at the points of multiplicities. In the following, we will assume that A is a matrix of (possibly) complex valued first order symbols, continuous with respect to t, i.e., A(t,x,ξ)(CS1)m×m.

We will now develop a parameter dependent extension of the Schur triangularisation procedure and we will describe it step by step. Then we will give an example for it for the systems of low sizes, namely, for m=2 and m=3.

In the case of m=2 the construction below was introduced in [27] and now we give its general version for systems of any size.

Normal forms of matrices depending on several parameters have a long history and are notoriously involved; for some remarks and related works, we refer the reader to [14, 24, 25, 45].

First step or Schur step

The first step in our triangularisation follows the construction in the constant case except that we will not get a unitary transformation matrix. For this reason we talk of a Schur step. Throughout this paper ei denotes the i-th vector of the standard basis of Rn with an appropriate dimension n.

Proposition 1

(Schur step) Let the m×m matrix valued symbol A(t,x,ξ)(CS1)m×m, have a real eigenvalue λCS1 and a corresponding eigenvector h(CS1)m such that there exists j{1,,m} with

h(t,x,ξ)|ej0(t,x,ξ)[0,T]×Rn×{|ξ|M}, 29

for a sufficiently large M>0. Then there exist an m×m matrix valued symbol T(t,x,ξ)(CS0)m×m, invertible for (t,x,ξ)[0,T]×Rn×{|ξ|M} with T-1(CS0)m×m, and an (m-1)×(m-1) matrix valued symbol E(t,x,ξ)(CS1)m×m, such that

T-1(t,x,ξ)A(t,x,ξ)T(t,x,ξ)=λa12a1m0E(t,x,ξ)0

for all (t,x,ξ)[0,T]×Rn×{|ξ|M}.

Proof

First let us note that we can assume that j=1 in (29). If that is not the case, we can exchange the rows 1 and j as well as columns 1 and j to move the jth component of the eigenvector to the first component.

We define the rescaled eigenvector μ componentwise by

μi(t,x,ξ)=h(t,x,ξ)|eih(t,x,ξ)|e1i=1,,m.

Now we set

T(t,x,ξ)=μ100μ2Im-1μm.

Since μ11 it follows that

T-1(t,x,ξ)=μ100-μ2Im-1-μm,

where Im-1 is the (m-1)×(m-1) identity matrix. By direct computations we get

AT=j=1ma1jμjA(2)A(m)j=1mamjμj=λμ1A(2)A(m)λμm,

where we used that

j=1maijμj=λμi,i=1,,m, 30

and denoted the ith column of A by A(i). The equations in (30) are given by the eigenvalue equation Aμ=λμ. Further, from μ11 we obtain

T-1AT=λμ12a12μ1a1mμ1-μ2μ1λ+μ2λE-μmμ1λ+μmλ=λa12a1m0E0, 31

which concludes the proof. Note that by construction the matrix E has entries in CS1 which depend on A. In particular its eigenvalues are the eigenvalues of A excluding λ (counted as many times as they occur).

Applying Proposition 1 repeatedly for m-2 times to E, we obtain a full Schur transformation of A, that is a full reduction to an upper triangular form. In the next subsection we describe this iteration in detail. This triangularisation procedure is summarised in Theorem 6 where sufficient conditions on the eigenvectors of A are given.

The triangularisation procedure

The reduction to an upper triangular form or the Schur transformation of A is possible under certain conditions on its eigenvectors. More precisely, let

h1(t,x,ξ),,hm-1(t,x,ξ)(CS0)m

be m-1 eigenvectors of A(t,x,ξ)=[aij(t,x,ξ)]i,j=1m, aijCS1, corresponding to the eigenvalues λ1(t,x,ξ), , λm-1(t,x,ξ)CS1. To formulate the sufficient conditions for the existence of such Schur transformation, we introduce a set of auxiliary vectors h(i), i=1,,m-1, which depend only on hi and the previous vectors h(j)CS0, j=1,,i-1. When i=1 we set h(1)=h1.

As in Proposition 1 we begin by assuming

h(1)(t,x,ξ)|e10 32

for (t,x,ξ)[0,T]×Rn×{|ξ|M}.

Remark 5

As noted in the proof of Proposition 1, we could have that

h(1)(t,x,ξ)|ej0

for another arbitrary j{1,,m}. Then, we could transform the matrix A(t,x,ξ) by a constant permutation matrix P such that P-1h(1) is eigenvector of P-1AP corresponding to λ1 which satisfies P-1h(1)(t,x,ξ)|e10. For this reason we state (32) with h(1) and e1.

Step 1
By Proposition 1 there exists a matrix T1 such that
T1-1AT1=λ1a12a1m0Em-10.
The matrix T1 is given by
T1=ω1e2em,ω1=ω11ω1mT
with
ω1j=h(1)(t,x,ξ)|ejh(1)(t,x,ξ)|e1.
In the sequel we make use of the projector Πk:RmRm-k, 0km-1, defined by
Πkx1xm=xk+1xm.
Note that Π0 is the identity map Im:RmRm.
Step 2
Since h2 is an eigenvector of A with eigenvalue λ2 we get that T1-1h2 is an eigenvector of T1-1AT1 with eigenvalue λ2 as well. By the structure of T1-1AT1 we easily see that h(2):=Π1T1-1h2 is an eigenvector of Em-1, corresponding to λ2. Arguing as in Remark 5 we assume that
Π1T1-1h2|e10(t,x,ξ)[0,T]×Rn×{|ξ|M}, 33
to be able to apply Proposition 1 to Em-1. We get that there exists an (m-1)×(m-1) matrix T~2 such that T~2-1Em-1T~2 is of form
λ20Em-20,
where in the first row the first row of Em-1 appears. Thus, setting
T2=1000T~20,
we obtain
T2-1T1-1AT1T2=λ10λ200Em-200. 34
Note that in (34) we write explicitly only the entries most relevant to our triangularisation. To compute the matrix T~2, we set
ω2=ω22ω2mT,
where
ω2j(t,x,ξ):=h(2)(t,x,ξ)|ejh(2)(t,x,ξ)|e1,j=2,,m,
and then
T~2=ω2e2em-1.
It is clear that T2 has the same structure as T1, i.e., it is defined via a rescaled eigenvector as the first column and an identity matrix (Im-1 for T1 and Im-2 for T2).
Step K
By iterating the method k-1 times we can find k-1 matrices T1,T2,,Tk-1 of size m×m such that
Tk-1-1··T1-1AT1··Tk-1=λ1000λk-1000Em-k+1000,
where Em-k+1 is a (m-k+1)×(m-k+1) matrix and the equality is true on [0,T]×Rn×{|ξ|M}. Since hk is an eigenvector of A corresponding to λk, the vector
Tk-1-1Tk-2-1··T1-1hk
is an eigenvector of
Tk-1-1Tk-2-1··T1-1AT1T2··Tk-1
and
h(k):=Πk-1Tk-1-1Tk-2-1··T1-1hk(CS0)m-k+1
an eigenvector of Em-k+1 corresponding to λk. Thus, to satisfy the assumptions of Proposition 1 and keeping in mind Remark 5, we require that
h(k)(t,x,ξ)|e10(t,x,ξ)[0,T]×Rn×{|ξ|M}. 35
It follows that there exists an (m-k+1)×(m-k+1) transformation matrix T~k such that T~k-1T~1-1AT~1T~k is of the form
λk0Em-k0.
and set
Tk=Ik-100T~k.
The matrix T~k is defined by
T~k=ωke2em-k+1,ωk=ωkkωkmT,
where
ωkj=h(k)(t,x,ξ)|ejh(k)(t,x,ξ)|e1,j=k,,m.
Stepm-1
This is the last step as E2 is a 2×2 matrix. We have that
h(m-1)=Πm-2Tm-2-1··T1-1hm-1(CS0)2
is an eigenvector of E2 corresponding to λm-1 and that T~m-1 exists as before if
h(m-1)(t,x,ξ)|e10(t,x,ξ)[0,T]×Rn×{|ξ|M}. 36
The matrix T~m-1 is given by
T~m-1=ωm-1e2=ωm-1,m-10ωm-1,m1,
where
ωm-1,j=h(m-1)(t,x,ξ)|ejh(m-1)(t,x,ξ)|e1,j=m-1,m,
and then
Tm-1=Im-200T~m-1.

We are now ready to state Theorem 6 which summarises the triangularisation procedure explained above. For the convenience of the reader we recall the notations introduced so far:

  • h1,,hm-1 are the eigenvectors of the matrix A corresponding to the eigenvalues λ1,,λm-1.

  • h(1)=h1 and
    h(i)=Πi-1Ti-1-1Ti-2-1··T1-1hi(CS0)m-k+1, 37
    for i=2,,m-1.
  • the matrices Tk are inductively defined as follows: T0=Im and
    Tk=Ik-100T~k,T~k=ωke2em-k,eiRm-k,
    where
    ωkj=h(k)(t,x,ξ)|ejh(k)(t,x,ξ)|e1,j=k,,m.

Finally, we note that h(k) depends only on Tk-1, , T1 and, thus, only on the eigenvectors h(k-1), , h(1).

Summarising, we can formulate a more precise version of Theorem 2.

Theorem 6

(Schur Decomposition) Let A(t,x,ξ)(CS1)m×m be a matrix with eigenvalues λ1,,λmCS1, and let h1,,hm-1(CS0)m be the corresponding eigenvectors. Suppose that for e1Rm-i+1 the condition

h(i)(t,x,ξ)|e10,(t,x,ξ)[0,T]×Rn×Rn 38

holds for all i=1,,m-1, with the notation explained above. Then, there exists a matrix-valued symbol T(t,x,ξ)(CS0)m×m, invertible for (t,x,ξ)[0,T]×Rn×{|ξ|M}, T-1(t,x,ξ)(CS0)m×m, such that

T-1(t,x,ξ)A(t,x,ξ)T(t,x,ξ)=Λ(t,x,ξ)+N(t,x,ξ)

for all (t,x,ξ)[0,T]×Rn×{|ξ|M}, where

Λ(t,x,ξ)=diag(λ1(t,x,ξ),λ2(t,x,ξ),,λm(t,x,ξ))

and

N(t,x,ξ)=0N12(t,x,ξ)N13(t,x,ξ)N1m(t,x,ξ)00N23(t,x,ξ)N2m(t,x,ξ)000Nm-1m(t,x,ξ)0000,

and N is a nilpotent matrix with entries in CS1. Furthermore, the matrix symbol T is given by

T(t,x,ξ)=T1T2··Tm-1,

with the notation explained above.

Remark 6

Taking into account Remark 5, let us stress that condition (38) is not restrictive as it can be replaced by the following: suppose that there exist m-1 numbers ji{1,,m-i+1}, i=1,m-1, such that for all i=1,,m-1

h(i)(t,x,ξ)|eji0(t,x,ξ)[0,T]×Rn×{|ξ|M} 39

holds.

Remark 7

If A(t,x,ξ) has complex symbols (as allowed in Theorem 1, see also Remark 1) and real eigenvalues, the eigenvalues of the Schur transformed system clearly remain real. The upper triangular entries may still be complex valued symbols.

Remark 8

Theorem 6 is quite general in the sense that the functions aij could be complex-valued. In this paper, we are concerned with hyperbolic matrices, i.e. we assume that the eigenvalues λ1,,λm are real. We stress that the Schur transform does not change the hyperbolicity of the matrix as the eigenvalues of T-1AT are also λ1,,λm.

Remark 9

For our applications in this and future work it is important that the transform T in Theorem 6 keeps the regularity of the original matrix A, i.e. that the elements of the Schur transform T-1AT are in the same class as the elements of A. Here, we stated everything with CS1 and CS0 as that is the regularity considered in this paper. Note that one could replace C with Ck or C and find a matrix T such that the transformed matrix T-1AT inherits the same regularity with respect to t. In addition, one could also drop the regularity in t to L and the triangularisation procedure would still work preserving the boundedness in t through every step.

For the sake of simplicity and the reader’s convenience, in the next subsections we analyse Theorem  6 in the special cases of m=2 and m=3.

The case m=2

We now formulate Theorem 6 in the special case m=2. In this way we recover the formulation given in [27].

Theorem 7

([27, Theorem 7.1]) Suppose that A(t,x,ξ)(CS1)2×2 admits eigenvalues λj(t,x,ξ)CS1, j=1,2, and an eigenvector h(t,x,ξ)(CS0)2 satisfying

h(t,x,ξ)|ej0,(t,x,ξ)[0,T]×Rn×{|ξ|M}, 40

for j=1 or j=2. Then, we can find a 2×2 matrix valued symbol T(t,x,ξ)(CS0)2×2, invertible for {|ξ|M} with T-1(t,x,ξ)(CS0)2×2, such that

T-1(t,x,ξ)A(t,x,ξ)T(t,x,ξ)=λ1(t,x,ξ)a12(t,x,ξ)0λ2(t,x,ξ)

for all (t,x,ξ)[0,T]×Rn×{|ξ|M}.

Proof

For 2×2 matrices the triangularisation procedure described in the previous subsection can stop at Step 1. By Remark 5, we may assume that (40) holds for the eigenvector h corresponding to λ1 and for j=1. We set h=h1 and h(1)=h1. The vector

ω1=ω11(t,x,ξ)ω12(t,x,ξ),ω1j(t,x,ξ)=h(1)(t,x,ξ)|ejh(1)(t,x,ξ)|e1,

belongs to CS0 and is an eigenvector of A associated to λ1. We then set

T1(t,x,ξ)=ω1e2=ω11(t,x,ξ)0ω12(t,x,ξ)1.

With that, we obtain

A(t,x,ξ)T1(t,x,ξ)=a11ω11+a12ω12a12a21ω11+a22ω12a22

and finally, with

T1-1(t,x,ξ)=ω11(t,x,ξ)0-ω12(t,x,ξ)1,

we obtain

T1-1(t,x,ξ)A(t,x,ξ)T1(t,x,ξ)=a11ω112+a12ω12ω11a12ω11-a11ω12ω11-a12ω122+a21ω11+a22ω12-ω12a12+a22

By construction, we have

a11ω11+a12ω12=λ1ω11,a21ω11+a22ω12=λ1ω12,

and ω1=1. This yields a11ω11+a12ω12=λ1ω11=λ1 and

-a11ω12ω11-a12ω122+a21ω11+a22ω12=-ω12(a11ω11+a12ω12)+a21ω11+a22ω12=-λ1ω12+λ1ω12=0..

Using a11+a22=λ1+λ2, we obtain

-ω12a12+a22=-ω12a12+a22+a11ω11-a11ω11=λ2.

Thus, we get that

T1-1(t,x,ξ)A(t,x,ξ)T1(t,x,ξ)=λ1(t,x,ξ)a12(t,x,ξ)0λ2(t,x,ξ)

for (t,x,ξ)[0,T]×Rn×{|ξ|M}. This concludes the proof.

Example

  • (i)
    By direct computations we can easily see that if h1=[h11h12]T=e1 then the matrix A is automatically in the upper triangular form. Indeed,
    a21h11+a22h12=λ1h12
    implies a21=0. A typical example (already discussed in [27]) is the Jordan block matrix
    A=0100,
    where λ1=0 is an eigenvalue with eigenvector h1=e1.
  • (ii)
    Condition (40) is trivially fulfilled when detA0 and A is of the form
    aa-a-a,
    for a=a(t,x,ξ). Indeed, also in this case one can take 0 as an eigenvalue with eigenvector h1=[11]T.

The case m=3

With the notation introduced in Sect. 3.2, we assume that the 3×3 matrix A(t,x,ξ)(CS1)3×3 admits three eigenvalues λi(t,x,ξ)CS1, i=1,2,3, and two corresponding eigenvectors hi(t,x,ξ)(CS0)3, i=1,2. Then, we set h(1):=h1 and, as in Remark 6 we suppose that there is a j1{1,2,3} with

h(1)(t,x,ξ)|ej10(t,x,ξ)[0,T]×Rn×{|ξ|M}.

Thus, we can set

ω1j(t,x,ξ)=h(1)(t,x,ξ)|ejh(1)(t,x,ξ)|ej1.

Now, we rearrange the matrix A such that the first component of ω1 becomes identically equal to 1. Then, with j2,j3{1,2,3}\{j1}, we can write

T1-1=ω1e2e3-1=ω1j100-ω1j210-ω1j301=100-ω1j210-ω1j301,

which leads to

T1-1h2=ω1j100-ω1j210-ω1j301h2j1h2j2h2j3=h2j1-ω1j2h2j1+h2j2-ω1j3h2j1+h2j3.

We then get

h(2)=Π1T1-1h2=-ω1j2h2j1+h2j2-ω1j3h2j1+h2j3

and the condition (38) that there exists j{1,2} such that

h(2)(t,x,ξ)|ej0(t,x,ξ)[0,T]×Rn×{|ξ|M}

translates to: either

-ω1j2h2j1+h2j20h2j2h1j1-h1j2h2j10 41

or

-ω1j3h2j1+h2j30h2j3h1j1-h1j3h2j10 42

holds. Thus, assuming that (41) holds, the matrix T~2 is given by

ω210ω211=10-ω1j3h2j1+h2j3-ω1j2h2j1+h2j21,ω2j=h(2)(t,x,ξ)|ejh(2)(t,x,ξ)|ej2,j=1,2,

and the matrix T2 by

1000ω2100ω221=1000100-ω1j3h2j1+h2j3-ω1j2h2j1+h2j21.

Thus, we obtain

T(t,x,ξ)=T1T2=100ω1j210ω1j3011000100-ω1j3h2j1+h2j3-ω1j2h2j1+h2j21. 43

If we have (42) instead of (41), then we would need a permutation matrix

Pj2j3=100001010

in (43), i.e.

T(t,x,ξ)=T1(t,x,ξ)Pj2j3T2(t,x,ξ)

and

T2(t,x,ξ)=1000100-ω1j2h2j1+h2j2-ω1j3h2j1+h2j31.

See also Remark 6.

Thus, we can state

Theorem 8

Suppose that A(t,x,ξ)(CS1)3×3 admits three eigenvalues λiCS1, i=1,2,3, and two corresponding eigenvectors hi(t,x,ξ)(CS1)3, i=1,2. Suppose that there exists a j1{1,2,3} such that

h1j1(t,x,ξ)0(t,x,ξ)[0,T]×Rn×{|ξ|M}. 44

Further suppose that there exists j2{1,2,3}\{j1} such that

h2j2h1j1-h1j2h2j10(t,x,ξ)[0,T]×Rn×{|ξ|M}. 45

Then, there exists a matrix-valued symbol T(t,x,ξ)(CS0)3×3, invertible for all (t,x,ξ)[0,T]×Rn×{|ξ|M} with T-1(t,x,ξ)(CS0)3×3, such that

T-1(t,x,ξ)A(t,x,ξ)T(t,x,ξ)=Λ(t,x,ξ)+N(t,x,ξ)

holds for all (t,x,ξ)[0,T]×Rn×{|ξ|M}, where Λ(t,x,ξ)=diag(λ1,λ2,λ3) and

N(t,x,ξ)=0N13(t,x,ξ)N13(t,x,ξ)00N23(t,x,ξ)000.

We end this subsection by discussing some examples of 3×3 matrices fulfilling the assumptions above on their eigenvalues.

Examples

  • (i)
    If the matrix A has eigenvectors
    h1=101andh2=110
    then conditions (44) and (45) are easily fulfilled with j1=1 and j2=2. Indeed, h11=1 and
    h22h11-h12h21=h22h11=1.
    More in general to satisfy (44) and (45) it would be enough to have two eigenvectors
    h1=h11h12h13andh2=h21h22h23
    with h110, h220 and h12=0.
  • (ii)
    A matrix with eigenvectors
    h1=101andh2=110
    has a special form. Indeed, for λ1 and λ2 eigenvalues corresponding to h1 and h2, respectively, by using the eigenvector equations we obtain
    a13=λ1-a11,a23=-a21,a33=λ1-a31,
    and
    a12=λ2-a11,a22=λ2-a21,a32=-a31.
    Hence
    A=a11λ2-a11λ1-a11a21λ2-a21-a21a31-a31λ1-a31.

Footnotes

Michael Ruzhansky was supported in parts by EPSRC Grant EP/R003025/1 and by the Leverhulme Grant RPG-2017-151. No new data was collected or generated during the course of research.

Inspired by our colleague and friend Todor Gramchev (1956–2015).

Contributor Information

Claudia Garetto, Email: c.garetto@lboro.ac.uk.

Christian Jäh, Email: c.jaeh@lboro.ac.uk.

Michael Ruzhansky, Email: m.ruzhansky@imperial.ac.uk.

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